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CompVis/

stable-diffusion-v1-4

Stable Diffusion is a latent text-to-image diffusion model capable of generating photo-realistic images given any text input.

Public
CompVis/stable-diffusion-v1-4 cover image

Input

text prompt

Please upload an image file

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Settings

negative text prompt. (Default: empty)

Num Images

number of images to generate (Default: 1, 1 ≤ num_images ≤ 4)

Num Inference Steps

number of denoising steps (Default: 25, 1 ≤ num_inference_steps ≤ 512)

Guidance Scale

classifier-free guidance, higher means follow prompt more closely (Default: 7.5, 1 ≤ guidance_scale ≤ 20)

Strength

how much to follow the input image. 1 means ignore the image, 0 means follow the image exactly (Default: 0.8, 0 ≤ strength ≤ 1)

image width in px 13

image height in px 13

random seed, empty means random (Default: empty, 0 ≤ seed)

use compel library for weighted prompts 2

coefficient applied to lora weights (Default: 1)

Output

generated image #0
Model Information